InvokeAI/docs/other/README-CompViz.md

---
title: CompViz-Readme
---

# _README from [CompViz/stable-diffusion](https://github.com/CompVis/stable-diffusion)_

_Stable Diffusion was made possible thanks to a collaboration with
[Stability AI](https://stability.ai/) and [Runway](https://runwayml.com/) and
builds upon our previous work:_

[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
[Robin Rombach](https://github.com/rromb)\*,
[Andreas Blattmann](https://github.com/ablattmann)\*,
[Dominik Lorenz](https://github.com/qp-qp)\,
[Patrick Esser](https://github.com/pesser),
[Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>

## **CVPR '22 Oral**

which is available on [GitHub](https://github.com/CompVis/latent-diffusion). PDF
at [arXiv](https://arxiv.org/abs/2112.10752). Please also visit our
[Project page](https://ommer-lab.com/research/latent-diffusion-models/).

![txt2img-stable2](../assets/stable-samples/txt2img/merged-0006.png)
[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
model. Thanks to a generous compute donation from
[Stability AI](https://stability.ai/) and support from
[LAION](https://laion.ai/), we were able to train a Latent Diffusion Model on
512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/)
database. Similar to Google's [Imagen](https://arxiv.org/abs/2205.11487), this
model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text
prompts. With its 860M UNet and 123M text encoder, the model is relatively
lightweight and runs on a GPU with at least 10GB VRAM. See
[this section](#stable-diffusion-v1) below and the
[model card](https://huggingface.co/CompVis/stable-diffusion).

## Requirements

A suitable [conda](https://conda.io/) environment named `ldm` can be created and
activated with:

```
conda env create
conda activate ldm
```

Note that the first line may be abbreviated `conda env create`, since conda will
look for `environment.yml` by default.

You can also update an existing
[latent diffusion](https://github.com/CompVis/latent-diffusion) environment by
running

```bash
conda install pytorch torchvision -c pytorch
pip install transformers==4.19.2
pip install -e .
```

## Stable Diffusion v1

Stable Diffusion v1 refers to a specific configuration of the model architecture
that uses a downsampling-factor 8 autoencoder with an 860M UNet and CLIP
ViT-L/14 text encoder for the diffusion model. The model was pretrained on
256x256 images and then finetuned on 512x512 images.

\*Note: Stable Diffusion v1 is a general text-to-image diffusion model and
therefore mirrors biases and (mis-)conceptions that are present in its training
data. Details on the training procedure and data, as well as the intended use of
the model can be found in the corresponding
[model card](https://huggingface.co/CompVis/stable-diffusion). Research into the
safe deployment of general text-to-image models is an ongoing effort. To prevent
misuse and harm, we currently provide access to the checkpoints only for
[academic research purposes upon request](https://stability.ai/academia-access-form).
**This is an experiment in safe and community-driven publication of a capable
and general text-to-image model. We are working on a public release with a more
permissive license that also incorporates ethical considerations.\***

[Request access to Stable Diffusion v1 checkpoints for academic research](https://stability.ai/academia-access-form)

### Weights

We currently provide three checkpoints, `sd-v1-1.ckpt`, `sd-v1-2.ckpt` and
`sd-v1-3.ckpt`, which were trained as follows,

- `sd-v1-1.ckpt`: 237k steps at resolution `256x256` on
  [laion2B-en](https://huggingface.co/datasets/laion/laion2B-en). 194k steps at
  resolution `512x512` on
  [laion-high-resolution](https://huggingface.co/datasets/laion/laion-high-resolution)
  (170M examples from LAION-5B with resolution `>= 1024x1024`).
- `sd-v1-2.ckpt`: Resumed from `sd-v1-1.ckpt`. 515k steps at resolution
  `512x512` on "laion-improved-aesthetics" (a subset of laion2B-en, filtered to
  images with an original size `>= 512x512`, estimated aesthetics score `> 5.0`,
  and an estimated watermark probability `< 0.5`. The watermark estimate is from
  the LAION-5B metadata, the aesthetics score is estimated using an
  [improved aesthetics estimator](https://github.com/christophschuhmann/improved-aesthetic-predictor)).
- `sd-v1-3.ckpt`: Resumed from `sd-v1-2.ckpt`. 195k steps at resolution
  `512x512` on "laion-improved-aesthetics" and 10\% dropping of the
  text-conditioning to improve
  [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).

Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling steps show the relative improvements of
the checkpoints: ![sd evaluation results](../assets/v1-variants-scores.jpg)

### Text-to-Image with Stable Diffusion

![txt2img-stable2](../assets/stable-samples/txt2img/merged-0005.png)
![txt2img-stable2](../assets/stable-samples/txt2img/merged-0007.png)

Stable Diffusion is a latent diffusion model conditioned on the (non-pooled)
text embeddings of a CLIP ViT-L/14 text encoder.

#### Sampling Script

After [obtaining the weights](#weights), link them

```
mkdir -p models/ldm/stable-diffusion-v1/
ln -s <path/to/model.ckpt> models/ldm/stable-diffusion-v1/model.ckpt
```

and sample with

```
python scripts/txt2img.py --prompt "a photograph of an astronaut riding a horse" --plms
```

By default, this uses a guidance scale of `--scale 7.5`,
[Katherine Crowson's implementation](https://github.com/CompVis/latent-diffusion/pull/51)
of the [PLMS](https://arxiv.org/abs/2202.09778) sampler, and renders images of
size 512x512 (which it was trained on) in 50 steps. All supported arguments are
listed below (type `python scripts/txt2img.py --help`).

```commandline
usage: txt2img.py [-h] [--prompt [PROMPT]] [--outdir [OUTDIR]] [--skip_grid] [--skip_save] [--ddim_steps DDIM_STEPS] [--plms] [--laion400m] [--fixed_code] [--ddim_eta DDIM_ETA] [--n_iter N_ITER] [--H H] [--W W] [--C C] [--f F] [--n_samples N_SAMPLES] [--n_rows N_ROWS]
                  [--scale SCALE] [--from-file FROM_FILE] [--config CONFIG] [--ckpt CKPT] [--seed SEED] [--precision {full,autocast}]

optional arguments:
  -h, --help            show this help message and exit
  --prompt [PROMPT]     the prompt to render
  --outdir [OUTDIR]     dir to write results to
  --skip_grid           do not save a grid, only individual samples. Helpful when evaluating lots of samples
  --skip_save           do not save individual samples. For speed measurements.
  --ddim_steps DDIM_STEPS
                        number of ddim sampling steps
  --plms                use plms sampling
  --laion400m           uses the LAION400M model
  --fixed_code          if enabled, uses the same starting code across samples
  --ddim_eta DDIM_ETA   ddim eta (eta=0.0 corresponds to deterministic sampling
  --n_iter N_ITER       sample this often
  --H H                 image height, in pixel space
  --W W                 image width, in pixel space
  --C C                 latent channels
  --f F                 downsampling factor
  --n_samples N_SAMPLES
                        how many samples to produce for each given prompt. A.k.a. batch size
                        (note that the seeds for each image in the batch will be unavailable)
  --n_rows N_ROWS       rows in the grid (default: n_samples)
  --scale SCALE         unconditional guidance scale: eps = eps(x, empty) + scale * (eps(x, cond) - eps(x, empty))
  --from-file FROM_FILE
                        if specified, load prompts from this file
  --config CONFIG       path to config which constructs model
  --ckpt CKPT           path to checkpoint of model
  --seed SEED           the seed (for reproducible sampling)
  --precision {full,autocast}
                        evaluate at this precision

```

Note: The inference config for all v1 versions is designed to be used with
EMA-only checkpoints. For this reason `use_ema=False` is set in the
configuration, otherwise the code will try to switch from non-EMA to EMA
weights. If you want to examine the effect of EMA vs no EMA, we provide "full"
checkpoints which contain both types of weights. For these, `use_ema=False` will
load and use the non-EMA weights.

#### Diffusers Integration

Another way to download and sample Stable Diffusion is by using the
[diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers)

```py
# make sure you're logged in with `huggingface-cli login`
from torch import autocast
from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler

pipe = StableDiffusionPipeline.from_pretrained(
	"CompVis/stable-diffusion-v1-3-diffusers",
	use_auth_token=True
)

prompt = "a photo of an astronaut riding a horse on mars"
with autocast("cuda"):
    image = pipe(prompt)["sample"][0]

image.save("astronaut_rides_horse.png")
```

### Image Modification with Stable Diffusion

By using a diffusion-denoising mechanism as first proposed by
[SDEdit](https://arxiv.org/abs/2108.01073), the model can be used for different
tasks such as text-guided image-to-image translation and upscaling. Similar to
the txt2img sampling script, we provide a script to perform image modification
with Stable Diffusion.

The following describes an example where a rough sketch made in
[Pinta](https://www.pinta-project.com/) is converted into a detailed artwork.

```
python scripts/img2img.py --prompt "A fantasy landscape, trending on artstation" --init-img <path-to-img.jpg> --strength 0.8
```

Here, strength is a value between 0.0 and 1.0, that controls the amount of noise
that is added to the input image. Values that approach 1.0 allow for lots of
variations but will also produce images that are not semantically consistent
with the input. See the following example.

**Input**

![sketch-in](../assets/stable-samples/img2img/sketch-mountains-input.jpg)

**Outputs**

![out3](../assets/stable-samples/img2img/mountains-3.png)
![out2](../assets/stable-samples/img2img/mountains-2.png)

This procedure can, for example, also be used to upscale samples from the base
model.

## Comments

- Our codebase for the diffusion models builds heavily on
  [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion) and
  [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch).
  Thanks for open-sourcing!

- The implementation of the transformer encoder is from
  [x-transformers](https://github.com/lucidrains/x-transformers) by
  [lucidrains](https://github.com/lucidrains?tab=repositories).

## BibTeX

```
@misc{rombach2021highresolution,
      title={High-Resolution Image Synthesis with Latent Diffusion Models},
      author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
      year={2021},
      eprint={2112.10752},
      archivePrefix={arXiv},
      primaryClass={cs.CV}
}

```