InvokeAI/README.md

# Stable Diffusion Dream Script

This is a fork of CompVis/stable-diffusion, the wonderful open source
text-to-image generator.

The original has been modified in several ways:

## Interactive command-line interface similar to the Discord bot

The *dream.py* script, located in scripts/dream.py, 
provides an interactive interface to image generation similar to
the "dream mothership" bot that Stable AI provided on its Discord
server. Unlike the txt2img.py and img2img.py scripts provided in the
original CompViz/stable-diffusion source code repository, the
time-consuming initialization of the AI model
initialization only happens once. After that image generation 
from the command-line interface is very fast.

The script uses the readline library to allow for in-line editing,
command history (up and down arrows), autocompletion, and more. To help
keep track of which prompts generated which images, the script writes a
log file of image names and prompts to the selected output directory.
In addition, as of version 1.02, it also writes the prompt into the PNG
file's metadata where it can be retrieved using scripts/images2prompt.py

The script is confirmed to work on Linux and Windows systems. It should
work on MacOSX as well, but this is not confirmed. Note that this script
runs from the command-line (CMD or Terminal window), and does not have a GUI.

~~~~
(ldm) ~/stable-diffusion$ python3 ./scripts/dream.py
* Initializing, be patient...
Loading model from models/ldm/text2img-large/model.ckpt
LatentDiffusion: Running in eps-prediction mode
DiffusionWrapper has 872.30 M params.
making attention of type 'vanilla' with 512 in_channels
Working with z of shape (1, 4, 32, 32) = 4096 dimensions.
making attention of type 'vanilla' with 512 in_channels
Loading Bert tokenizer from "models/bert"
setting sampler to plms

* Initialization done! Awaiting your command...
dream> ashley judd riding a camel -n2 -s150
Outputs:
   outputs/img-samples/00009.png: "ashley judd riding a camel" -n2 -s150 -S 416354203
   outputs/img-samples/00010.png: "ashley judd riding a camel" -n2 -s150 -S 1362479620

dream> "there's a fly in my soup" -n6 -g
    outputs/img-samples/00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
    seeds for individual rows: [2685670268, 1216708065, 2335773498, 822223658, 714542046, 3395302430]
dream> q

# this shows how to retrieve the prompt stored in the saved image's metadata
(ldm) ~/stable-diffusion$ python3 ./scripts/images2prompt.py outputs/img_samples/*.png
00009.png: "ashley judd riding a camel" -s150 -S 416354203
00010.png: "ashley judd riding a camel" -s150 -S 1362479620
00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
~~~~

The dream> prompt's arguments are pretty much identical to those used
in the Discord bot, except you don't need to type "!dream" (it doesn't
hurt if you do). A significant change is that creation of individual
images is now the default unless --grid (-g) is given. For backward
compatibility, the -i switch is recognized.  For command-line help
type -h (or --help) at the dream> prompt.

The script itself also recognizes a series of command-line switches
that will change important global defaults, such as the directory for
image outputs and the location of the model weight files.

## Image-to-Image

This script also provides an img2img feature that lets you seed your
creations with a drawing or photo. This is a really cool feature that tells
stable diffusion to build the prompt on top of the image you provide, preserving
the original's basic shape and layout. To use it, provide the --init_img 
option as shown here:

~~~~
dream> "waterfall and rainbow" --init_img=./init-images/crude_drawing.png --strength=0.5 -s100 -n4
~~~~

The --init_img (-I) option gives the path to the seed picture. --strength (-f) controls how much
the original will be modified, ranging from 0.0 (keep the original intact), to 1.0 (ignore the original
completely). The default is 0.75, and ranges from 0.25-0.75 give interesting results.

## Weighted Prompts

You may weight different sections of the prompt to tell the sampler to attach different levels of
priority to them, by adding :(number) to the end of the section you wish to up- or downweight.
For example consider this prompt:

~~~~
    tabby cat:0.25 white duck:0.75 hybrid
~~~~

This will tell the sampler to invest 25% of its effort on the tabby
cat aspect of the image and 75% on the white duck aspect
(surprisingly, this example actually works). The prompt weights can
use any combination of integers and floating point numbers, and they
do not need to add up to 1. A practical example of using this type of
weighting is described here:
https://www.reddit.com/r/StableDiffusion/comments/wvb7q7/using_prompt_weights_to_tweak_an_image_with/

## Changes

* v1.06 (23 August 2022)
   * Added weighted prompt support contributed by [xraxra](https://github.com/xraxra)
   * Example of using weighted prompts to tweak a demonic figure contributed by [bmaltais](https://github.com/bmaltais)

* v1.05 (22 August 2022 - after the drop)
   * Filenames now use the following formats:
       000010.95183149.png      -- Two files produced by the same command (e.g. -n2),
       000010.26742632.png      -- distinguished by a different seed.

       000011.455191342.01.png  -- Two files produced by the same command using
       000011.455191342.02.png  -- a batch size>1 (e.g. -b2). They have the same seed.

       000011.4160627868.grid#1-4.png  -- a grid of four images (-g); the whole grid can
                                          be regenerated with the indicated key

    * It should no longer be possible for one image to overwrite another
    * You can use the "cd" and "pwd" commands at the dream> prompt to set and retrieve
      the path of the output directory.
     
* v1.04 (22 August 2022 - after the drop)
   * Updated README to reflect installation of the released weights.
   * Suppressed very noisy and inconsequential warning when loading the frozen CLIP
   tokenizer.

* v1.03 (22 August 2022)
   * The original txt2img and img2img scripts from the CompViz repository have been moved into
     a subfolder named "orig_scripts", to reduce confusion.
   
* v1.02 (21 August 2022)
    * A copy of the prompt and all of its switches and options is now stored in the corresponding
    image in a tEXt metadata field named "Dream". You can read the prompt using scripts/images2prompt.py,
    or an image editor that allows you to explore the full metadata.
        **Please run "conda env update -f environment.yaml" to load the k_lms dependencies!!**

* v1.01 (21 August 2022)
    * added k_lms sampling. 
      **Please run "conda env update -f environment.yaml" to load the k_lms dependencies!!**
    * use half precision arithmetic by default, resulting in faster execution and lower memory requirements
    Pass argument --full_precision to dream.py to get slower but more accurate image generation


## Installation

There are separate installation walkthroughs for [Linux/Mac](#linuxmac) and [Windows](#windows).

### Linux/Mac

1. You will need to install the following prerequisites if they are not already available. Use your
operating system's preferred installer
* Python (version 3.8.5 recommended; higher may work)
* git

2. Install the Python Anaconda environment manager using pip3.
```
~$ pip3 install anaconda
```
After installing anaconda, you should log out of your system and log back in. If the installation
worked, your command prompt will be prefixed by the name of the current anaconda environment, "(base)".

3. Copy the stable-diffusion source code from GitHub:
```
(base) ~$ git clone https://github.com/lstein/stable-diffusion.git
```
This will create stable-diffusion folder where you will follow the rest of the steps.

4. Enter the newly-created stable-diffusion folder. From this step forward make sure that you are working in the stable-diffusion directory!
```
(base) ~$ cd stable-diffusion
(base) ~/stable-diffusion$
```
5. Use anaconda to copy necessary python packages, create a new python environment named "ldm",
and activate the environment. 
```
(base) ~/stable-diffusion$ conda env create -f environment.yaml
(base) ~/stable-diffusion$ conda activate ldm
(ldm) ~/stable-diffusion$
```
After these steps, your command prompt will be prefixed by "(ldm)" as shown above.

6. Load a couple of small machine-learning models required by stable diffusion:
```
(ldm) ~/stable-diffusion$ python3 scripts/preload_models.py
```

Note that this step is necessary because I modified the original
just-in-time model loading scheme to allow the script to work on GPU
machines that are not internet connected. See [Workaround for machines with limited internet connectivity](#workaround-for-machines-with-limited-internet-connectivity)

7. Now you need to install the weights for the stable diffusion model.

For running with the released weights, you will first need to set up an acount with Hugging Face (https://huggingface.co).
Use your credentials to log in, and then point your browser at https://huggingface.co/CompVis/stable-diffusion-v-1-4-original.
You may be asked to sign a license agreement at this point.

Click on "Files and versions" near the top of the page, and then click on the file named "sd-v1-4.ckpt". You'll be taken
to a page that prompts you to click the "download" link. Save the file somewhere safe on your local machine.

Now run the following commands from within the stable-diffusion directory. This will create a symbolic
link from the stable-diffusion model.ckpt file, to the true location of the sd-v1-4.ckpt file.
   
```
(ldm) ~/stable-diffusion$ mkdir -p models/ldm/stable-diffusion-v1
(ldm) ~/stable-diffusion$ ln -sf /path/to/sd-v1-4.ckpt models/ldm/stable-diffusion-v1/model.ckpt
```

8. Start generating images!
```
# for the pre-release weights use the -l or --liaon400m switch
(ldm) ~/stable-diffusion$ python3 scripts/dream.py -l

# for the post-release weights do not use the switch
(ldm) ~/stable-diffusion$ python3 scripts/dream.py

# for additional configuration switches and arguments, use -h or --help
(ldm) ~/stable-diffusion$ python3 scripts/dream.py -h
```
9. Subsequently, to relaunch the script, be sure to run "conda activate ldm" (step 5, second command), enter the "stable-diffusion" 
directory, and then launch the dream script (step 8). If you forget to activate the ldm environment, the script will fail with multiple ModuleNotFound errors.

#### Updating to newer versions of the script

This distribution is changing rapidly. If you used the "git clone" method (step 5) to download the stable-diffusion directory, then to update to the latest and greatest version, launch the Anaconda window, enter "stable-diffusion", and type:
```
(ldm) ~/stable-diffusion$ git pull
```
This will bring your local copy into sync with the remote one.

### Windows

1. Install Python version 3.8.5 from here: https://www.python.org/downloads/windows/
   (note that several users have reported that later versions do not work properly)

2. Install Anaconda3 (miniconda3 version) from here: https://docs.anaconda.com/anaconda/install/windows/

3. Install Git from here: https://git-scm.com/download/win

4. Launch Anaconda from the Windows Start menu. This will bring up a command window. Type all the remaining commands in this window.

5. Run the command:
```
git clone https://github.com/lstein/stable-diffusion.git
```
This will create stable-diffusion folder where you will follow the rest of the steps.

6. Enter the newly-created stable-diffusion folder. From this step forward make sure that you are working in the stable-diffusion directory!
```
cd stable-diffusion
```

7. Run the following two commands:
```
conda env create -f environment.yaml    (step 7a)
conda activate ldm                      (step 7b)
```
This will install all python requirements and activate the "ldm" environment which sets PATH and other environment variables properly.

8. Run the command:
```
python scripts\preload_models.py
```

This installs several machine learning models that stable diffusion
requires. (Note that this step is required. I created it because some people
are using GPU systems that are behind a firewall and the models can't be
downloaded just-in-time)

9. Now you need to install the weights for the big stable diffusion model.

For running with the released weights, you will first need to set up
an acount with Hugging Face (https://huggingface.co).  Use your
credentials to log in, and then point your browser at
https://huggingface.co/CompVis/stable-diffusion-v-1-4-original.  You
may be asked to sign a license agreement at this point.

Click on "Files and versions" near the top of the page, and then click
on the file named "sd-v1-4.ckpt". You'll be taken to a page that
prompts you to click the "download" link. Now save the file somewhere
safe on your local machine.  The weight file is >4 GB in size, so
downloading may take a while.

Now run the following commands from **within the stable-diffusion
directory** to copy the weights file to the right place:
   
```
mkdir -p models/ldm/stable-diffusion-v1
copy C:\path\to\sd-v1-4.ckpt models\ldm\stable-diffusion-v1\model.ckpt
```
Please replace "C:\path\to\sd-v1.4.ckpt" with the correct path to wherever
you stashed this file. If you prefer not to copy or move the .ckpt file, 
you may instead create a shortcut to it from within
"models\ldm\stable-diffusion-v1\".

10. Start generating images!
```
# for the pre-release weights
python scripts\dream.py -l

# for the post-release weights
python scripts\dream.py
```
11. Subsequently, to relaunch the script, first activate the Anaconda command window (step 4), enter the stable-diffusion directory (step 6, "cd \path\to\stable-diffusion"), run "conda activate ldm" (step 7b), and then launch the dream script (step 10).

#### Updating to newer versions of the script

This distribution is changing rapidly. If you used the "git clone" method (step 5) to download the stable-diffusion directory, then to update to the latest and greatest version, launch the Anaconda window, enter "stable-diffusion", and type:
```
git pull
```
This will bring your local copy into sync with the remote one.

## Simplified API for text to image generation

For programmers who wish to incorporate stable-diffusion into other
products, this repository includes a simplified API for text to image generation, which
lets you create images from a prompt in just three lines of code:

~~~~
from ldm.simplet2i import T2I
model   = T2I()
outputs = model.txt2img("a unicorn in manhattan")
~~~~

Outputs is a list of lists in the format [[filename1,seed1],[filename2,seed2]...]
Please see ldm/simplet2i.py for more information.


## Workaround for machines with limited internet connectivity

My development machine is a GPU node in a high-performance compute
cluster which has no connection to the internet. During model
initialization, stable-diffusion tries to download the Bert tokenizer
and a file needed by the kornia library.  This obviously didn't work
for me.

To work around this, I have modified ldm/modules/encoders/modules.py
to look for locally cached Bert files rather than attempting to
download them. For this to work, you must run
"scripts/preload_models.py" once from an internet-connected machine
prior to running the code on an isolated one. This assumes that both
machines share a common network-mounted filesystem with a common
.cache directory.

~~~~
(ldm) ~/stable-diffusion$ python3 ./scripts/preload_models.py
preloading bert tokenizer...
Downloading: 100%|██████████████████████████████████| 28.0/28.0 [00:00<00:00, 49.3kB/s]
Downloading: 100%|██████████████████████████████████| 226k/226k [00:00<00:00, 2.79MB/s]
Downloading: 100%|██████████████████████████████████| 455k/455k [00:00<00:00, 4.36MB/s]
Downloading: 100%|██████████████████████████████████| 570/570 [00:00<00:00, 477kB/s]
...success
preloading kornia requirements...
Downloading: "https://github.com/DagnyT/hardnet/raw/master/pretrained/train_liberty_with_aug/checkpoint_liberty_with_aug.pth" to /u/lstein/.cache/torch/hub/checkpoints/checkpoint_liberty_with_aug.pth
100%|███████████████████████████████████████████████| 5.10M/5.10M [00:00<00:00, 101MB/s]
...success
~~~~

If you don't need this change and want to download the files just in
time, copy over the file ldm/modules/encoders/modules.py from the
CompVis/stable-diffusion repository. Or you can run preload_models.py
on the target machine.

## Support

For support,
please use this repository's GitHub Issues tracking service. Feel free
to send me an email if you use and like the script.

*Original Author:* Lincoln D. Stein <lincoln.stein@gmail.com>

*Contributions by:* [Peter Kowalczyk](https://github.com/slix), [Henry Harrison](https://github.com/hwharrison), [xraxra](https://github.com/xraxra), and [bmaltais](https://github.com/bmaltais)

# Original README from CompViz/stable-diffusion
*Stable Diffusion was made possible thanks to a collaboration with [Stability AI](https://stability.ai/) and [Runway](https://runwayml.com/) and builds upon our previous work:*

[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
[Robin Rombach](https://github.com/rromb)\*,
[Andreas Blattmann](https://github.com/ablattmann)\*,
[Dominik Lorenz](https://github.com/qp-qp)\,
[Patrick Esser](https://github.com/pesser),
[Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>

**CVPR '22 Oral**

which is available on [GitHub](https://github.com/CompVis/latent-diffusion). PDF at [arXiv](https://arxiv.org/abs/2112.10752). Please also visit our [Project page](https://ommer-lab.com/research/latent-diffusion-models/).

![txt2img-stable2](assets/stable-samples/txt2img/merged-0006.png)
[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
model.
Thanks to a generous compute donation from [Stability AI](https://stability.ai/) and support from [LAION](https://laion.ai/), we were able to train a Latent Diffusion Model on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database. 
Similar to Google's [Imagen](https://arxiv.org/abs/2205.11487), 
this model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts.
With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 10GB VRAM.
See [this section](#stable-diffusion-v1) below and the [model card](https://huggingface.co/CompVis/stable-diffusion).

  
## Requirements

A suitable [conda](https://conda.io/) environment named `ldm` can be created
and activated with:

```
conda env create -f environment.yaml
conda activate ldm
```

You can also update an existing [latent diffusion](https://github.com/CompVis/latent-diffusion) environment by running

```
conda install pytorch torchvision -c pytorch
pip install transformers==4.19.2
pip install -e .
```

## Stable Diffusion v1

Stable Diffusion v1 refers to a specific configuration of the model
architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet
and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and 
then finetuned on 512x512 images.

*Note: Stable Diffusion v1 is a general text-to-image diffusion model and therefore mirrors biases and (mis-)conceptions that are present
in its training data. 
Details on the training procedure and data, as well as the intended use of the model can be found in the corresponding [model card](https://huggingface.co/CompVis/stable-diffusion).
Research into the safe deployment of general text-to-image models is an ongoing effort. To prevent misuse and harm, we currently provide access to the checkpoints only for [academic research purposes upon request](https://stability.ai/academia-access-form).
**This is an experiment in safe and community-driven publication of a capable and general text-to-image model. We are working on a public release with a more permissive license that also incorporates ethical considerations.***

[Request access to Stable Diffusion v1 checkpoints for academic research](https://stability.ai/academia-access-form) 

### Weights

We currently provide three checkpoints, `sd-v1-1.ckpt`, `sd-v1-2.ckpt` and `sd-v1-3.ckpt`,
which were trained as follows,

- `sd-v1-1.ckpt`: 237k steps at resolution `256x256` on [laion2B-en](https://huggingface.co/datasets/laion/laion2B-en).
  194k steps at resolution `512x512` on [laion-high-resolution](https://huggingface.co/datasets/laion/laion-high-resolution) (170M examples from LAION-5B with resolution `>= 1024x1024`).
- `sd-v1-2.ckpt`: Resumed from `sd-v1-1.ckpt`.
  515k steps at resolution `512x512` on "laion-improved-aesthetics" (a subset of laion2B-en,
filtered to images with an original size `>= 512x512`, estimated aesthetics score `> 5.0`, and an estimated watermark probability `< 0.5`. The watermark estimate is from the LAION-5B metadata, the aesthetics score is estimated using an [improved aesthetics estimator](https://github.com/christophschuhmann/improved-aesthetic-predictor)).
- `sd-v1-3.ckpt`: Resumed from `sd-v1-2.ckpt`. 195k steps at resolution `512x512` on "laion-improved-aesthetics" and 10\% dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).

Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling
steps show the relative improvements of the checkpoints:
![sd evaluation results](assets/v1-variants-scores.jpg)



### Text-to-Image with Stable Diffusion
![txt2img-stable2](assets/stable-samples/txt2img/merged-0005.png)
![txt2img-stable2](assets/stable-samples/txt2img/merged-0007.png)

Stable Diffusion is a latent diffusion model conditioned on the (non-pooled) text embeddings of a CLIP ViT-L/14 text encoder.


#### Sampling Script

After [obtaining the weights](#weights), link them
```
mkdir -p models/ldm/stable-diffusion-v1/
ln -s <path/to/model.ckpt> models/ldm/stable-diffusion-v1/model.ckpt 
```
and sample with
```
python scripts/txt2img.py --prompt "a photograph of an astronaut riding a horse" --plms 
```

By default, this uses a guidance scale of `--scale 7.5`, [Katherine Crowson's implementation](https://github.com/CompVis/latent-diffusion/pull/51) of the [PLMS](https://arxiv.org/abs/2202.09778) sampler, 
and renders images of size 512x512 (which it was trained on) in 50 steps. All supported arguments are listed below (type `python scripts/txt2img.py --help`).

```commandline
usage: txt2img.py [-h] [--prompt [PROMPT]] [--outdir [OUTDIR]] [--skip_grid] [--skip_save] [--ddim_steps DDIM_STEPS] [--plms] [--laion400m] [--fixed_code] [--ddim_eta DDIM_ETA] [--n_iter N_ITER] [--H H] [--W W] [--C C] [--f F] [--n_samples N_SAMPLES] [--n_rows N_ROWS]
                  [--scale SCALE] [--from-file FROM_FILE] [--config CONFIG] [--ckpt CKPT] [--seed SEED] [--precision {full,autocast}]

optional arguments:
  -h, --help            show this help message and exit
  --prompt [PROMPT]     the prompt to render
  --outdir [OUTDIR]     dir to write results to
  --skip_grid           do not save a grid, only individual samples. Helpful when evaluating lots of samples
  --skip_save           do not save individual samples. For speed measurements.
  --ddim_steps DDIM_STEPS
                        number of ddim sampling steps
  --plms                use plms sampling
  --laion400m           uses the LAION400M model
  --fixed_code          if enabled, uses the same starting code across samples
  --ddim_eta DDIM_ETA   ddim eta (eta=0.0 corresponds to deterministic sampling
  --n_iter N_ITER       sample this often
  --H H                 image height, in pixel space
  --W W                 image width, in pixel space
  --C C                 latent channels
  --f F                 downsampling factor
  --n_samples N_SAMPLES
                        how many samples to produce for each given prompt. A.k.a. batch size
  --n_rows N_ROWS       rows in the grid (default: n_samples)
  --scale SCALE         unconditional guidance scale: eps = eps(x, empty) + scale * (eps(x, cond) - eps(x, empty))
  --from-file FROM_FILE
                        if specified, load prompts from this file
  --config CONFIG       path to config which constructs model
  --ckpt CKPT           path to checkpoint of model
  --seed SEED           the seed (for reproducible sampling)
  --precision {full,autocast}
                        evaluate at this precision

```
Note: The inference config for all v1 versions is designed to be used with EMA-only checkpoints. 
For this reason `use_ema=False` is set in the configuration, otherwise the code will try to switch from
non-EMA to EMA weights. If you want to examine the effect of EMA vs no EMA, we provide "full" checkpoints
which contain both types of weights. For these, `use_ema=False` will load and use the non-EMA weights.


#### Diffusers Integration

Another way to download and sample Stable Diffusion is by using the [diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers)
```py
# make sure you're logged in with `huggingface-cli login`
from torch import autocast
from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler

pipe = StableDiffusionPipeline.from_pretrained(
	"CompVis/stable-diffusion-v1-3-diffusers", 
	use_auth_token=True
)

prompt = "a photo of an astronaut riding a horse on mars"
with autocast("cuda"):
    image = pipe(prompt)["sample"][0]  
    
image.save("astronaut_rides_horse.png")
```



### Image Modification with Stable Diffusion

By using a diffusion-denoising mechanism as first proposed by [SDEdit](https://arxiv.org/abs/2108.01073), the model can be used for different 
tasks such as text-guided image-to-image translation and upscaling. Similar to the txt2img sampling script, 
we provide a script to perform image modification with Stable Diffusion.  

The following describes an example where a rough sketch made in [Pinta](https://www.pinta-project.com/) is converted into a detailed artwork.
```
python scripts/img2img.py --prompt "A fantasy landscape, trending on artstation" --init-img <path-to-img.jpg> --strength 0.8
```
Here, strength is a value between 0.0 and 1.0, that controls the amount of noise that is added to the input image. 
Values that approach 1.0 allow for lots of variations but will also produce images that are not semantically consistent with the input. See the following example.

**Input**

![sketch-in](assets/stable-samples/img2img/sketch-mountains-input.jpg)

**Outputs**

![out3](assets/stable-samples/img2img/mountains-3.png)
![out2](assets/stable-samples/img2img/mountains-2.png)

This procedure can, for example, also be used to upscale samples from the base model.


## Comments 

- Our codebase for the diffusion models builds heavily on [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion)
and [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch). 
Thanks for open-sourcing!

- The implementation of the transformer encoder is from [x-transformers](https://github.com/lucidrains/x-transformers) by [lucidrains](https://github.com/lucidrains?tab=repositories). 


## BibTeX

```
@misc{rombach2021highresolution,
      title={High-Resolution Image Synthesis with Latent Diffusion Models}, 
      author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
      year={2021},
      eprint={2112.10752},
      archivePrefix={arXiv},
      primaryClass={cs.CV}
}

```