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* update IMG2IMG.md * update INPAINTING.md * update WEBUIHOTKEYS.md * more doc updates (mostly fix formatting): - OUTPAINTING.md - POSTPROCESS.md - PROMPTS.md - VARIATIONS.md - WEB.md - WEBUIHOTKEYS.md
236 lines
10 KiB
Markdown
236 lines
10 KiB
Markdown
---
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title: Image-to-Image
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---
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# :material-image-multiple: Image-to-Image
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## `img2img`
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This script also provides an `img2img` feature that lets you seed your creations
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with an initial drawing or photo. This is a really cool feature that tells
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stable diffusion to build the prompt on top of the image you provide, preserving
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the original's basic shape and layout. To use it, provide the `--init_img`
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option as shown here:
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```commandline
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tree on a hill with a river, nature photograph, national geographic -I./test-pictures/tree-and-river-sketch.png -f 0.85
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```
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This will take the original image shown here:
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<figure markdown>
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![](https://user-images.githubusercontent.com/50542132/193946000-c42a96d8-5a74-4f8a-b4c3-5213e6cadcce.png)
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</figure>
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and generate a new image based on it as shown here:
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<figure markdown>
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![](https://user-images.githubusercontent.com/111189/194135515-53d4c060-e994-4016-8121-7c685e281ac9.png)
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</figure>
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The `--init_img` (`-I`) option gives the path to the seed picture. `--strength`
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(`-f`) controls how much the original will be modified, ranging from `0.0` (keep
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the original intact), to `1.0` (ignore the original completely). The default is
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`0.75`, and ranges from `0.25-0.90` give interesting results. Other relevant
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options include `-C` (classification free guidance scale), and `-s` (steps).
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Unlike `txt2img`, adding steps will continuously change the resulting image and
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it will not converge.
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You may also pass a `-v<variation_amount>` option to generate `-n<iterations>`
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count variants on the original image. This is done by passing the first
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generated image back into img2img the requested number of times. It generates
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interesting variants.
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Note that the prompt makes a big difference. For example, this slight variation
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on the prompt produces a very different image:
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<figure markdown>
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![](https://user-images.githubusercontent.com/111189/194135220-16b62181-b60c-4248-8989-4834a8fd7fbd.png)
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<caption markdown>photograph of a tree on a hill with a river</caption>
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</figure>
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!!! tip
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When designing prompts, think about how the images scraped from the internet were captioned. Very few photographs will
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be labeled "photograph" or "photorealistic." They will, however, be captioned with the publication, photographer, camera
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model, or film settings.
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If the initial image contains transparent regions, then Stable Diffusion will
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only draw within the transparent regions, a process called
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[`inpainting`](./INPAINTING.md#creating-transparent-regions-for-inpainting).
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However, for this to work correctly, the color information underneath the
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transparent needs to be preserved, not erased.
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!!! warning
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**IMPORTANT ISSUE** `img2img` does not work properly on initial images smaller
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than 512x512. Please scale your image to at least 512x512 before using it.
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Larger images are not a problem, but may run out of VRAM on your GPU card. To
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fix this, use the --fit option, which downscales the initial image to fit within
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the box specified by width x height:
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```
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tree on a hill with a river, national geographic -I./test-pictures/big-sketch.png -H512 -W512 --fit
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```
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## How does it actually work, though?
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The main difference between `img2img` and `prompt2img` is the starting point.
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While `prompt2img` always starts with pure gaussian noise and progressively
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refines it over the requested number of steps, `img2img` skips some of these
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earlier steps (how many it skips is indirectly controlled by the `--strength`
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parameter), and uses instead your initial image mixed with gaussian noise as the
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starting image.
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**Let's start** by thinking about vanilla `prompt2img`, just generating an image
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from a prompt. If the step count is 10, then the "latent space" (Stable
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Diffusion's internal representation of the image) for the prompt "fire" with
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seed `1592514025` develops something like this:
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```commandline
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invoke> "fire" -s10 -W384 -H384 -S1592514025
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```
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<figure markdown>
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![latent steps](../assets/img2img/000019.steps.png)
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</figure>
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Put simply: starting from a frame of fuzz/static, SD finds details in each frame
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that it thinks look like "fire" and brings them a little bit more into focus,
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gradually scrubbing out the fuzz until a clear image remains.
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**When you use `img2img`** some of the earlier steps are cut, and instead an
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initial image of your choice is used. But because of how the maths behind Stable
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Diffusion works, this image needs to be mixed with just the right amount of
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noise (fuzz/static) for where it is being inserted. This is where the strength
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parameter comes in. Depending on the set strength, your image will be inserted
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into the sequence at the appropriate point, with just the right amount of noise.
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### A concrete example
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I want SD to draw a fire based on this hand-drawn image:
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<figure markdown>
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![drawing of a fireplace](../assets/img2img/fire-drawing.png)
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</figure>
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Let's only do 10 steps, to make it easier to see what's happening. If strength
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is `0.7`, this is what the internal steps the algorithm has to take will look
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like:
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<figure markdown>
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![gravity32](../assets/img2img/000032.steps.gravity.png)
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</figure>
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With strength `0.4`, the steps look more like this:
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<figure markdown>
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![gravity30](../assets/img2img/000030.steps.gravity.png)
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</figure>
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Notice how much more fuzzy the starting image is for strength `0.7` compared to
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`0.4`, and notice also how much longer the sequence is with `0.7`:
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| | strength = 0.7 | strength = 0.4 |
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| --------------------------- | ------------------------------------------------------------- | ------------------------------------------------------------- |
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| initial image that SD sees | ![](../assets/img2img/000032.step-0.png) | ![](../assets/img2img/000030.step-0.png) |
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| steps argument to `invoke>` | `-S10` | `-S10` |
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| steps actually taken | 7 | 4 |
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| latent space at each step | ![gravity32](../assets/img2img/000032.steps.gravity.png) | ![gravity30](../assets/img2img/000030.steps.gravity.png) |
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| output | ![000032.1592514025](../assets/img2img/000032.1592514025.png) | ![000030.1592514025](../assets/img2img/000030.1592514025.png) |
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Both of the outputs look kind of like what I was thinking of. With the strength
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higher, my input becomes more vague, _and_ Stable Diffusion has more steps to
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refine its output. But it's not really making what I want, which is a picture of
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cheery open fire. With the strength lower, my input is more clear, _but_ Stable
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Diffusion has less chance to refine itself, so the result ends up inheriting all
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the problems of my bad drawing.
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If you want to try this out yourself, all of these are using a seed of
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`1592514025` with a width/height of `384`, step count `10`, the default sampler
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(`k_lms`), and the single-word prompt `"fire"`:
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```commandline
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invoke> "fire" -s10 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png --strength 0.7
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```
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The code for rendering intermediates is on my (damian0815's) branch
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[document-img2img](https://github.com/damian0815/InvokeAI/tree/document-img2img) -
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run `invoke.py` and check your `outputs/img-samples/intermediates` folder while
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generating an image.
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### Compensating for the reduced step count
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After putting this guide together I was curious to see how the difference would
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be if I increased the step count to compensate, so that SD could have the same
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amount of steps to develop the image regardless of the strength. So I ran the
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generation again using the same seed, but this time adapting the step count to
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give each generation 20 steps.
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Here's strength `0.4` (note step count `50`, which is `20 ÷ 0.4` to make sure SD
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does `20` steps from my image):
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```commandline
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invoke> "fire" -s50 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png -f 0.4
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```
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<figure markdown>
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![000035.1592514025](../assets/img2img/000035.1592514025.png)
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</figure>
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and here is strength `0.7` (note step count `30`, which is roughly `20 ÷ 0.7` to
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make sure SD does `20` steps from my image):
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```commandline
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invoke> "fire" -s30 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png -f 0.7
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```
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<figure markdown>
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![000046.1592514025](../assets/img2img/000046.1592514025.png)
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</figure>
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In both cases the image is nice and clean and "finished", but because at
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strength `0.7` Stable Diffusion has been give so much more freedom to improve on
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my badly-drawn flames, they've come out looking much better. You can really see
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the difference when looking at the latent steps. There's more noise on the first
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image with strength `0.7`:
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<figure markdown>
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![gravity46](../assets/img2img/000046.steps.gravity.png)
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</figure>
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than there is for strength `0.4`:
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<figure markdown>
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![gravity35](../assets/img2img/000035.steps.gravity.png)
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</figure>
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and that extra noise gives the algorithm more choices when it is evaluating how
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to denoise any particular pixel in the image.
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Unfortunately, it seems that `img2img` is very sensitive to the step count.
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Here's strength `0.7` with a step count of `29` (SD did 19 steps from my image):
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<figure markdown>
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![gravity45](../assets/img2img/000045.1592514025.png)
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</figure>
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By comparing the latents we can sort of see that something got interpreted
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differently enough on the third or fourth step to lead to a rather different
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interpretation of the flames.
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<figure markdown>
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![gravity46](../assets/img2img/000046.steps.gravity.png)
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</figure>
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<figure markdown>
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![gravity45](../assets/img2img/000045.steps.gravity.png)
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</figure>
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This is the result of a difference in the de-noising "schedule" - basically the
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noise has to be cleaned by a certain degree each step or the model won't
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"converge" on the image properly (see
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[stable diffusion blog](https://huggingface.co/blog/stable_diffusion) for more
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about that). A different step count means a different schedule, which means
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things get interpreted slightly differently at every step.
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